Stable Diffusion
Stable Diffusion is a deep learning, text-to-image model released in 2022 based on diffusion techniques. It is considered to be a part of the ongoing AI spring.
An image generated by Stable Diffusion based on the text prompt "a photograph of an astronaut riding a horse" | |
| Original author(s) | Runway, CompVis, and Stability AI |
|---|---|
| Developer(s) | Stability AI |
| Initial release | August 22, 2022 |
| Stable release | SDXL 1.0 (model)
/ July 26, 2023 |
| Repository | |
| Written in | Python |
| Operating system | Any that support CUDA kernels |
| Type | Text-to-image model |
| License | Creative ML OpenRAIL-M |
| Website | stability |
It is primarily used to generate detailed images conditioned on text descriptions, though it can also be applied to other tasks such as inpainting, outpainting, and generating image-to-image translations guided by a text prompt. Its development involved researchers from the CompVis Group at Ludwig Maximilian University of Munich and Runway with a computational donation by Stability AI and training data from non-profit organizations.
Stable Diffusion is a latent diffusion model, a kind of deep generative artificial neural network. Its code and model weights have been open sourced, and it can run on most consumer hardware equipped with a modest GPU with at least 4 GB VRAM. This marked a departure from previous proprietary text-to-image models such as DALL-E and Midjourney which were accessible only via cloud services.